Tag: diffusion
Diffusion probabilistic models were originally proposed as a latent variable generative model inspired by non-equilibrium thermodynamics. The essential idea of diffusion models is to systematically perturb the structure in a data distribution through a forward diffusion process, and then recover the structure by learning a reverse diffusion process, resulting in a highly flexible and tractable generative model.- A Picture is Worth a Thousand Words Principled Recaptioning Improves Image Generation (28 Oct 2023)
This is my reading note for A Picture is Worth a Thousand Words: Principled Recaptioning Improves Image Generation. The papers found that the text data used to train text to image model is now quality, which is based alt text of images.it proposed to use an image caption model to generate high quality text for the images; then the diffusion model trained from this new text data show much better performance.
- Idea2Img Iterative Self-Refinement with GPT-4V(ision) for Automatic Image Design and Generation (14 Oct 2023)
This is my reading note for Idea2Img: Iterative Self-Refinement with GPT-4V(ision) for Automatic Image Design and Generation. This paper proposes a system on how to use GPT4V to generate images from idea by calling an image generation tool. Especially.it generates text prompt based on idea, given the images generated from the prompt, it ranks and selects the best image; it then generate a new promote to guide image generation process.
- Word-As-Image for Semantic Typography (13 Oct 2023)
This is my reading note for Word-As-Image for Semantic Typography. This paper utilized the differential rendering for vector graph to train a diffusion model to generate vector graph for a given text. Check my note for related paper in # VectorFusion Text-to-SVG by Abstracting Pixel-Based Diffusion Models
- VectorFusion Text-to-SVG by Abstracting Pixel-Based Diffusion Models (12 Oct 2023)
This is my reading note for VectorFusion: Text-to-SVG by Abstracting Pixel-Based Diffusion Models. This paper utilized the differential rendering for vector graph to train a diffusion model to generate vector graph for a given text.
- DiffusionDet Diffusion Model for Object Detection (06 Oct 2023)
This is my reading note for DiffusionDet: Diffusion Model for Object Detection. This paper formulates the object detection problem as a diffusion process: recover object bounding box from noisy estimation. The initial estimation could be from purely random Gaussian noise. One benefit of this method is that it could automatically handle different number of bounding boxes
- Raising the Cost of Malicious AI-Powered Image Editing (03 Oct 2023)
This is my reading note for Raising the Cost of Malicious AI-Powered Image Editing. This paper proposes a method to stop an image being edited by on diffusion model. The method is based on adverbial attack: learn a perturbation to the target image such that the model (encoder or diffusion) will generate noise or degraded image. However this method may not always work or may fall when the model changes.
- FreeU Free Lunch in Diffusion U-Net (20 Sep 2023)
This is my reading note for FreeU: Free Lunch in Diffusion U-Net. The paper analyzed the cause of artifact from diffusion model. The paper should that the backbone (U-Net) captures the global or low frequency information and skip connection capture the fine detail or high frequency information.it also shows that the high frequency information causes artifacts. As a results, this paper proposes increasing weight of half channel of U-Net and suppress the low frequency information from the skip connection
- 360 Reconstruction From a Single Image Using Space Carved Outpainting (19 Sep 2023)
This is my reading note for 360 Reconstruction From a Single Image Using Space Carved Outpainting. This paper proposes a method of 3D reconstruction from a single image. To the it represents the 3D object by NERF and iteratively update the NERF by rendering new view using Dream booth.
- Rerender A Video Zero-Shot Text-Guided Video-to-Video Translation (18 Sep 2023)
This is my reading note on Rerender A Video: Zero-Shot Text-Guided Video-to-Video Translation. The paper proposes a method to edit a video given style mentioned in prompt. The method performed diffusion to edit key frames and then propagate the edited key frames to other frames using optical flow. For key frame editing, several attention based constraint is applied to reserve details and consistency, including shape aware, style aware, pixel aware and fidelity aware.
- NExT-GPT Any-to-Any Multimodal LLM (16 Sep 2023)
This is my reading note for NExT-GPT: Any-to-Any Multimodal LLM. This paper proposes a multiple modality model which could takes multiple modalities as input and output in multiple modalities as well. The paper leverage existing large language model, multiple modality encoder image bind) and multiple modality diffusion model. To Amish the spice of those components, a simple linear projection is used for input and transformer to the output.
- PhotoVerse Tuning-Free Image Customization with Text-to-Image Diffusion Models (13 Sep 2023)
This is my reading note for PhotoVerse: Tuning-Free Image Customization with Text-to-Image Diffusion. This paper proposes a fine tune free personalized image edit method bases on diffusion. To this end it proposes dual branch to encode text and image feature. Lora is used to update the existing model. it also proposed to use a random fusion injection to condition the attention with image and text embedding. However the paper fails to describe why this random fusion injection is needed.
- MagiCapture High-Resolution Multi-Concept Portrait Customization (11 Sep 2023)
This is my reading note on MagiCapture High-Resolution Multi-Concept Portrait Customization. This paper proposes a diffusion method to apply a style to a specific face image. Both the style and face are given as images. To do this, this paper fine tune existing model with LORA given several new loss functions: one is face identity loss for the face region given a face recognition model; another one is background similarity for the style. The two loss are applied to the latent vector.
- InstructDiffusion A Generalist Modeling Interface for Vision Tasks (10 Sep 2023)
This is my reading note for InstructDiffusion: A Generalist Modeling Interface for Vision Tasks. This paper formulated many vision tasks like segmentation and key point detection as text guided image edit task, and thus can be modeled by diffusion based image edit model. To to that, this paper collects a dataset of different vision tasks, each item contains source image, vision task as text prompt and target image as vision results.
- InstaFlow One Step is Enough for High-Quality Diffusion-Based Text-to-Image Generation (09 Sep 2023)
This is my reading note on InstaFlow One Step is Enough for High-Quality Diffusion-Based Text-to-Image. This paper proposes a way to speed up diffusion based method, by achieving high fidelity with one step of diffusion. The key to this method is to use rectified how to straighten the probability flow from model to the final image. After that the model could be distilled to one step diffusion.
- Key-Locked Rank One Editing for Text-to-Image Personalization (07 Sep 2023)
This is my reading note on Key-Locked Rank One Editing for Text-to-Image Personalization. This paper proposes a personalized image generation method base on controlling attention module of the diffusion model. Especially key captures the layout of concept and value captures the identity of the new concept. A rank one update is applied to the attention weight to this purpose.
- DiffBIR Towards Blind Image Restoration with Generative Diffusion Prior (06 Sep 2023)
This is my reading note on DiffBIR Towards Blind Image Restoration with Generative Diffusion Prior. This paper proposes a two stage method for restore degraded images: stage 1 is trained neural network to recover image degradation; stage 2 is a pretrained diffusion model to restore the details in the image recovered from stage 1.
- DreamFusion Text-to-3D using 2D Diffusion (01 Sep 2023)
This is my reading note on DreamFusion: Text-to-3D using 2D Diffusion. This paper proposes a method (score distillation sampling or SDS) to distill a pre-trained text to image diffusion model to a 3D model. The 3D model, which is based on NERF, is trained per text prompt.
- DreamBooth Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation (31 Aug 2023)
This is my reading note on DreamBooth. Given as input just a few images of a subject, we fine-tune a pretrained text-to-image model (Imagen, although our method is not limited to a specific model) such that it learns to bind a unique identifier with that specific subject. Once the subject is embedded in the output domain of the model, the unique identifier can then be used to synthesize fully-novel photorealistic images of the subject contextualized in different scenes. By leveraging the semantic prior embedded in the model with a new autogenous class-specific prior preservation loss, our technique enables synthesizing the subject in diverse scenes, poses, views, and lighting conditions that do not appear in the reference images.
- TokenFlow Consistent Diffusion Features for Consistent Video Editing (29 Aug 2023)
This is my reading note on TokenFlow Consistent Diffusion Features for Consistent Video Editing, which is diffusion based on video editing method. This paper proposes a method to edit a video given text prompt. To do this, the paper relies on two things. First, it extracts bey lames from video and perform image on those key frames jointly. In addition, the paper found that the feature in diffusion has strong correspondence to the pixels. As a results it propose to propagate the features of edited key frames to other frames, accord to the correspondence in the original video.
- Diff-Instruct A Universal Approach for Transferring Knowledge From Pre-trained Diffusion Models (28 Aug 2023)
This is my reading note on Diff-Instruct: A Universal Approach for Transferring Knowledge From Pre-trained Diffusion Models. The paper explains the theory of using a pre-trained diffusion model to guide the training of a generator model.it shows that both DreamFusion and GAN are a special case of it: score distillation sampling (SDS) from DreamFusion uses Dirac distribution to represent the generator while GAN learns a discriminator to represents the distribution of data. To this end, it proposes IKL, which is tailored for DMs by calculating the integral of the KL divergence along a diffusion process (instead of a single step), which we show to be more robust in comparing distributions with misaligned supports.
- StableVideo Text-driven Consistency-aware Diffusion Video Editing (22 Aug 2023)
This is my reading note on StableVideo: Text-driven Consistency-aware Diffusion Video Editing. This paper proposes a video editing method based on diffusion. To ensure temporal consistency, the method utilizes neural atlas and inter frame interpolation. The neural atlas separate the videos into foreground and background plane. The lattes defines the mapping of pixel in frame to u v coordinate in atlas. For inter frame interpolation, the edited imago from diffusion is mapping to next frame via atlas, which is then use as initial to denote to the final contents of this frame.
- BLIP-Diffusion Pre-trained Subject Representation for Controllable Text-to-Image Generation and Editing (21 Aug 2023)
This is my reading note for BLIP-Diffusion: Pre-trained Subject Representation for Controllable Text-to-Image Generation and Editing. The paper proposes a method for generating an image with text prompt and target visual concept. To do that the paper trained blip model to align visual features with text prompt and then concatenate the visual embedding to the text prompt to generate the need. Code and models will be released at https://github.com/salesforce/LAVIS/tree/main/projects/blip-diffusion. Project page at https://dxli94.github.io/BLIP-Diffusion-website/.
- ProlificDreamer High-Fidelity and Diverse Text-to-3D Generation with Variational Score Distillation (20 Aug 2023)
This is my reading note on ProlificDreamer High-Fidelity and Diverse Text-to-3D Generation with Variational Score Distillation. This method proposes variational score sampling to replace score distillation sampling to improve the details of text to image or text to 3D models. Project page: https://ml.cs.tsinghua.edu.cn/prolificdreamer/
- Elucidating the Design Space of Diffusion-Based Generative Models (19 Aug 2023)
This is my reading note for Elucidating the Design Space of Diffusion-Based Generative Models. This paper checks the varying design of diffusion method and proposed a unify frame work to incorporate them. Finally the author proposes optimal choice of diffusion method under this frame work.
- Scalable Adaptive Computation for Iterative Generation (18 Aug 2023)
This is my reading note on Scalable Adaptive Computation for Iterative Generation The major innovation here is to map the input token to latents, which is shorter. The latents could be initialized from previous iterations (of diffusion process). As a result, the new method could achieve similar visual fidelity as regular diffusion method but with 1/10 of cost.
Introduction
- NeuralField-LDM Scene Generation with Hierarchical Latent Diffusion Models (17 Aug 2023)
This is my reading note on NeuralField-LDM Scene Generation with Hierarchical Latent Diffusion Models. It trains auto-encoder to project RGB images of scene with camera pose into the latent space (voxel-nerf). It uses three levels of latent to represent the scene and then uses hierarchical latent diffusion model to represent it.
- BLIP-2 Bootstrapping Language-Image Pre-training with Frozen Image Encoders and Large Language Models (04 Aug 2023)
This is my reading note for BLIP-2: Bootstrapping Language-Image Pre-training with Frozen Image Encoders and Large Language Models. The paper propose Q former to align the visual feature to text feature. Both visual feature and text feature are extracted from fixed models. Q former learned query and output the visual embeds to the text space.
- HyperDreamBooth HyperNetworks for Fast Personalization of Text-to-Image Models (27 Jul 2023)
This is my reading note for HyperDreamBooth: HyperNetworks for Fast Personalization of Text-to-Image Models. This paper improves DreamBooth by applying LORA to improve speed.
- Subject-Diffusion Open Domain Personalized Text-to-Image Generation without Test-time Fine-tuning (26 Jul 2023)
This is my reading note for Subject-Diffusion:Open Domain Personalized Text-to-Image Generation without Test-time Fine-tuning. This paper propose a diffusion method to generate images with given visual concepts and text prompt. Especially the paper is able to hand multiple visual concert jointly. To handle that, the paper detect the visual concepts from the input images, then the segmented images and bounding box are encoded feed into latent diffusion model. To enhance the consistency, the visual embedding is inserted into the text encode of the prompt.
- Make-An-Audio Text-To-Audio Generation with Prompt-Enhanced Diffusion Models (23 Jul 2023)
This is my reading note for Make-An-Audio: Text-To-Audio Generation with Prompt-Enhanced Diffusion Models. This paper proposes a diffusion model for audio, which uses an auto encoder to convert audio signal to a spectrum which could be natively handled by latent diffusion method.
- DreamBooth Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation (14 Jul 2023)
This is my reading note for DreamBooth: Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation. This paper proposes a personalized method for text to image based on diffusion. To achieve this, it firsts learn to align the visual content to be personalized to a rarely used text embedding, then this text embedding will be insert to the text to control the image generation.
- Aligning Text-to-Image Diffusion Models with Reward Backpropagation (10 Jul 2023)
This is my reading note for Aligning Text-to-Image Diffusion Models with Reward Backpropagation. This paper proposes a method how to train diffusion model for a given reward function in a memory efficient way, especially it utilities Lora and checkpoints . To avoid model collapse, it also proposes to randomly truncate number of steps.
- Blended Latent Diffusion (06 Jul 2023)
This is my reading note for Blended Latent Diffusion. The major innovation of the paper is to apply mask in latent space instead of image space to reduce boundary inconsistency, as the foreground is generated from the VAE but the background is not. in addition to handle the thin detail of mask got lost due to downs sample step, it dilate the mask first.
- Blended Latent Diffusion (05 Jul 2023)
This is my reading note for Blended Latent Diffusion. The major innovation of the paper is to apply mask in latent space instead of image space to reduce boundary inconsistency, as the foreground is generated from the VAE but the background is not. in addition to handle the thin detail of mask got lost due to downs sample step, it dilate the mask first.
- MeshDiffusion Score-based Generative 3D Mesh Modeling (02 Jul 2023)
This is my reading note for MeshDiffusion: Score-based Generative 3D Mesh Modeling. This paper represents the 3D mesh as a reformed tetrahedral which is defined on a regular 3D grid with 4 channel features: 3D positional deformation of the vertex and signed distance function values to define the surface.
- Scaling Autoregressive Multi-Modal Models Pretraining and Instruction Tuning (01 Jul 2023)
This is my reading note for Scaling Autoregressive Multi-Modal Models: Pretraining and Instruction Tuning. This paper proposes a method for text to image generation which is NOT based on diffusion. It utilizes auto-regressive model on tokens.
- Localizing and Editing Knowledge in Text-to-Image Generative Models (27 Jun 2023)
This is my reading note for Localizing and Editing Knowledge in Text-to-Image Generative Models. This paper studied how each component of diffusion model contribute to the final result: only that self attention layer of last tokens contribute to the final result. Then it proposes a simple method to perform image editing by modifying that layer.
- An Image is Worth Multiple Words Learning Object Level Concepts using Multi-Concept Prompt Learning (26 Jun 2023)
This is my reading note for An Image is Worth Multiple Words: Learning Object Level Concepts using Multi-Concept Prompt Learning. This paper proposes a method to learn embedding of multiple concepts for diffusion model, to this ends, it leverages masking in embed and contrast loss.
- DALL-E, DALL-E2 and StoryDALL-E (30 Sep 2022)
This my reading note on Zero-Shot Text-to-Image Generation (aka, DALL-E), its extension Hierarchical Text-Conditional Image Generation with CLIP Latents (aka, DALLE-2 or unCLIP) and StoryDALL-E: Adapting Pretrained Text-to-Image Transformers for Story Continuation. DALL-E is a transformer generating image given captions, by autoregressively modeling the text and image tokens as a single stream of data. StoryDALL-E extends DALL-E by generating a sequence of images for a sequence of caption to complete a story.
- DreamBooth Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation (28 Sep 2022)
This is my reading note on DreamBooth: Fine Tuning Text-to-Image Diffusion Models for Subject-Driven Generation. Given as input just a few (3~5) images of a subject, DreamBooth fine-tune a pretrained text-to-image model such that it learns to bind a unique identifier with that specific subject. Once the subject is embedded in the output domain of the model, the unique identifier can then be used to synthesize fully-novel photorealistic images of the subject contextualized in different scenes. By leveraging the semantic prior embedded in the model with a new autogenous class-specific prior preservation loss, DreamBooth enables synthesizing the subject in diverse scenes, poses, views, and lighting conditions that do not appear in the reference images. (check Figure 1 as an example)
- Recent Adavances of Diffusion Models (24 Sep 2022)
This is my 4th note in Diffusion models. For the previous notes, please refer to diffusion and stable diffusion. My contents are based on paper listed in Diffusion Explained and Diffusion Models: A Comprehensive Survey of Methods and Applications.
- unCLIP-Hierarchical Text-Conditional Image Generation with CLIP Latents (23 Sep 2022)
This is my reading note on Hierarchical Text-Conditional Image Generation with CLIP Latents. This paper proposes a two-stage model (unCLIP): a prior that generates a CLIP image embedding given a text caption, and a decoder that generates an image conditioned on the image embedding, for generating images from text.
- Stable Diffusion (23 Sep 2022)
This is my 2nd reading note on diffusion model, which will focus on the
stabe diffusion
, aka High-Resolution Image Synthesis with Latent Diffusion Models. By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. However, as mentioned in diffusion, DM sufferes high computational cost. The proposed Latent Diffusion Models (LDM) reduces the computational cost via latent space and introduces cross-attention to enable multi-modality conditioning. - Diffusion Model (22 Sep 2022)
This is my 1st reading note of on recent progress of difussion model. It is based on Diffusion Models: A Comprehensive Survey of Methods and Applications. Diffusion probabilistic models were originally proposed as a latent variable generative model inspired by non- equilibrium thermodynamics. The essential idea of diffusion models is to systematically perturb the structure in a data distribution through a forward diffusion process, and then recover the structure by learning a reverse diffusion process, resulting in a highly flexible and tractable generative model.